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Stable Diffusion (v1)

An open-weight text-to-image model released publicly by Stability AI on 22 August 2022, built on latent diffusion and small enough to run on consumer GPUs, which made local image generation widely accessible.

Status AvailableModality Text, ImageOpen weights Yes Era Large models

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RunwayStability AI

What it is

Stable Diffusion is an open-weight text-to-image model: it takes a written prompt and produces an image. Stability AI announced its public release on 22 August 2022, after a first-stage release to researchers announced on 10 August[1][4]. Because the weights were published openly, the model could be run on one's own hardware rather than only through a company's hosted service.

Where it came from

The method comes from latent diffusion, described in the paper "High-Resolution Image Synthesis with Latent Diffusion Models" by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser and Björn Ommer, posted to arXiv in December 2021 and published at CVPR 2022[3]. Hugging Face described the model as created by researchers and engineers from CompVis, Stability AI and LAION[5]; Stability's launch announcement also credited Runway and the Machine Vision & Learning group at LMU Munich[4].

How it works

A diffusion model learns to generate images by learning to reverse a noising process. During training, noise is added to images step by step and the network learns to remove it; at generation time the process starts from random noise and denoises, and an image appears. The latent-diffusion idea is to do this not on pixels but in the compressed latent space of an autoencoder: the autoencoder squeezes an image into a smaller representation, the diffusion process runs there, and a decoder turns the result back into a full image[3]. Working in that smaller space is what makes training and sampling affordable.

The text prompt enters through cross-attention. In version 1 the prompt is encoded by a frozen CLIP ViT-L/14 text encoder with 123 million parameters, and the result is fed into an 860 million-parameter UNet through cross-attention layers[2]. That is how words steer what the denoising process produces.

Training and released weights

The model was trained on a subset of the LAION-5B image-text database at 512x512 resolution[2]. Several checkpoints were published, sd-v1-1 through sd-v1-4, with v1.4 the recommended public weights[2]. Stability said training ran on its 4,000-A100 Ezra-1 cluster[4].

What shipped with it

The release came under the CreativeML OpenRAIL-M license, which allows commercial and non-commercial use together with a set of use restrictions[1]. The software package included an AI-based safety classifier, enabled by default, to filter unwanted outputs[1]. Stability also offered a hosted app, DreamStudio, and recommended it for more control and rapid generation[1].

Local use was meant to be within reach: Stability said the model runs on consumer GPUs with under 10 GB of VRAM[4], and that final memory usage at release should be about 6.9 GB[1].

Where it sits

Adoption was fast. By 17 October 2022, Stability said the model had been downloaded and licensed by more than 200,000 developers[6]. It also became a starting point for later models: the successor lines include Stable Diffusion 2.0 in November 2022 and SDXL, whose paper in July 2023 describes a UNet backbone three times larger along with a second text encoder[7].

CapabilitiesImage generation

TechniquesGenerative models

ApplicationsCreative tools

Ecosystem & regionOpen sourceEurope

Version history

5
  1. 2022
    • Aug 22, 2022 Stable Diffusion (v1) · This page

      An open-weight text-to-image model released publicly by Stability AI on 22 August 2022, built on latent diffusion and small enough to run on consumer GPUs, which made local image generation widely accessible.

    • Nov 24, 2022 Stable Diffusion 2

      Stable Diffusion 2 is an open-source text-to-image model family released by Stability AI on 24 November 2022, built on a new OpenCLIP text encoder and shipped with upscaling, depth-guided and inpainting models; version 2.1 followed weeks later.

  2. 2023
    • Jul 26, 2023 SDXL (Stable Diffusion XL)

      SDXL is Stability AI's open text-to-image diffusion model, announced in July 2023 with a 3.5B-parameter base model and a refiner stage, generating natively at 1024×1024. A distilled version, SDXL Turbo, cut generation down to a single step.

  3. 2024
    • Feb 22, 2024 Stable Diffusion 3

      Stable Diffusion 3, announced by Stability AI on 22 February 2024, is a suite of 800 million to 8 billion parameter models that pairs a diffusion transformer with flow matching.

    • Oct 22, 2024 Stable Diffusion 3.5

      A family of text-to-image models from Stability AI, announced in October 2024: the 8.1-billion-parameter Large, a 4-step Turbo variant and the 2.5-billion-parameter Medium, released under a community licence.

20251

  1. PolicyUK High Court rules on Getty Images v Stability AI

    On 4 November 2025 the High Court of England and Wales ruled in Getty Images v Stability AI, rejecting Getty's secondary copyright claim but finding very limited trademark infringement over watermarks in some outputs.

20231

  1. PolicyArtists sue Stability AI, Midjourney and DeviantArt over training data

    In January 2023 three artists filed a class-action copyright suit against Stability AI, Midjourney and DeviantArt, saying their image generators were trained on artists' work without permission — an early test of how copyright law applies to generative AI.

20222

  1. FundingStability AI raises $101 million

    On 17 October 2022 Stability AI announced $101 million in funding, led by Coatue, Lightspeed Venture Partners and O'Shaughnessy Ventures, to accelerate open AI models across several media.

  2. ReleaseStability AI publicly releases Stable Diffusion

    On 22 August 2022 Stability AI publicly released Stable Diffusion, a latent text-to-image diffusion model whose weights, code and model card were published on Hugging Face.

Developed by2

RunwayStability AI

Builds on1

CLIP

Succeeded by1

Stable Diffusion 2

Notes and sources

  1. Stable Diffusion Public Releasehttps://stability.ai/news-updates/stable-diffusion-public-release
  2. CompVis/stable-diffusion repositoryhttps://github.com/CompVis/stable-diffusion
  3. High-Resolution Image Synthesis with Latent Diffusion Modelshttps://arxiv.org/abs/2112.10752
  4. Stable Diffusion Launch Announcementhttps://stability.ai/news-updates/stable-diffusion-announcement
  5. Stable Diffusion with Diffusershttps://huggingface.co/blog/stable_diffusion
  6. Stability AI Announces $101 Million in Funding for Open-Source Artificial Intelligencehttps://stability.ai/news-updates/stability-ai-announces-101-million-in-funding-for-open-source-artificial-intelligence
  7. SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesishttps://arxiv.org/abs/2307.01952

Cite this page: Stable Diffusion (v1). TellWorld, https://www.tellworld.net/models/stable-diffusion. Updated 2026-11-08